SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis
Written for the Research library (claude-opus-5-5) from the paper text. Original — arxiv.org/pdf/2307.01952
Summary
Stable Diffusion is a widely used open latent diffusion model for text-to-image generation, but its earlier releases (SD 1.5, SD 2.1) lagged behind closed commercial systems in image quality. Stability AI researchers present SDXL, which enlarges the UNet about threefold and restructures it. Transformer blocks are concentrated at lower-resolution feature levels, and text conditioning comes from two encoders, CLIP ViT-L and OpenCLIP ViT-bigG, plus a pooled text embedding. Two lightweight extra inputs, the original image size and the crop offsets applied during data loading, let the model learn from small images instead of discarding them and reduce cut-off objects in outputs. The model is then finetuned on many aspect ratios. An optional refinement model applies SDEdit-style noising and denoising to the base latents to sharpen details. In user preference studies SDXL clearly beats earlier Stable Diffusion versions, the refiner adds further gains, and quality approaches that of black-box generators. Code and weights are openly released, which supports reproducibility and independent bias analysis.
Key points
- The SDXL UNet has 2.6B parameters (vs. 860M for SD 1.4/1.5 and 865M for SD 2.0/2.1), and its two text encoders add 817M more. It uses a transformer-block layout of [0, 2, 10], a 2048-dim context, and no 8× downsampling level.
- Without size-conditioning, 39% of the pretraining data would be dropped for having an edge shorter than 256 px.
- On class-conditional ImageNet at 512², size-conditioning gave the best scores: FID-5k 36.53 / IS-5k 215.34, compared with 39.76 / 211.50 (no conditioning) and 43.84 / 110.64 (training only on images ≥512 px).
- Generation runs in two stages. The base model produces 128 × 128 latents, which a high-resolution refiner sharing the same autoencoder then refines. Multi-aspect finetuning covers ratios from 0.25 to 4.0 at roughly 1024² pixels.
- On 10k COCO prompts, SDXL has the worst FID of the three compared models and only slightly better CLIP scores, despite strong human preference. The authors argue this shows the need for better automatic metrics for text-to-image models.
Cited by
- HunyuanVideo: A Systematic Framework For Large Video Generative Models
- The Dawn of LMMs: Preliminary Explorations with GPT-4V(ision)
- CogVideoX: Text-to-Video Diffusion Models with An Expert Transformer
Cite
@misc{podell2023sdxl,
title = {{SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis}},
author = {Podell, Dustin and English, Zion and Lacey, Kyle and Blattmann, Andreas and Dockhorn, Tim and Muller, Jonas and Penna, Joe and Rombach, Robin},
year = {2023},
doi = {10.48550/arxiv.2307.01952},
eprint = {2307.01952},
archivePrefix = {arXiv},
url = {https://doi.org/10.48550/arxiv.2307.01952}
}References
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